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bge-m3-GGUF

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**Model creator**: BAAI **Original model**: bge-m3 **GGUF quantization**: based on llama.cpp release 61408e7f For more details please refer to our github repo: https://github.com/FlagOpen/FlagEmbedding In this project, we introduce BGE-M3, which is distinguished for its versatility in Multi-Functionality, Multi-Linguality, and Multi-Granularity. - Multi-Functionality: It can simultaneously perform the three common retrieval functionalities of embedding model: dense retrieval, multi-vector retrieval, and sparse retrieval. - Multi-Linguality: It can support more than 100 working languages. - Multi-Granularity: It is able to process inputs of different granularities, spanning from short sentences to long documents of up to 8192 tokens.

Embeddings & similarity
Input Text→ Output Embedding
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**Model creator**: BAAI **Original model**: bge-reranker-v2-m3 **GGUF quantization**: based on llama.cpp release f4d2b **More details please refer to our Github: FlagEmbedding.** - Model List - Usage - Fine-tuning - Evaluation - Citation

Text analysis
Input Text→ Output Text
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基础档案已收录，详细介绍后续补充。可先查看厂商、分类与来源。

Category pending
Input Not verified→ Output Not verified
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**Model creator**: Runway ML **Original model**: stable-diffusion-v1-5 **GGUF quantization**: based on stable-diffusion.cpp ac54e that patched by llama-box. Stable Diffusion is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input. For more information about how Stable Diffusion functions, please have a look at 🤗's Stable Diffusion blog. The **Stable-Diffusion-v1-5** checkpoint was initialized with the weights of the Stable-Diffusion-v1-2 checkpoint and subsequently fine-tuned on 595k steps at resolution 512x512 on "laion-aesthetics v2 5+" and 10% dropping of the text-conditioning to improve classifier-free guidance sampling.

Image generation
Input Text→ Output Image
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**Model creator**: Stability AI **Original model**: stable-diffusion-xl-base-1.0 **GGUF quantization**: based on stable-diffusion.cpp ac54e that patched by llama-box. **VAE From**: madebyollin/sdxl-vae-fp16-fix. SDXL consists of an ensemble of experts pipeline for latent diffusion: In a first step, the base model is used to generate (noisy) latents, which are then further processed with a refinement model (available here: https://huggingface.co/stabilityai/stable-diffusion-xl-refiner-1.0/) specialized for the final denoising steps. Note that the base model can be used as a standalone module. Alternatively, we can use a two-stage pipeline as follows: First, the base model is used to generate latents of the desired output size. In the second step, we use a specialized high-resolution model and apply a technique called SDEdit (https://arxiv.org/abs/2108.01073, also known as "img2img") to the latents generated in the first step, using the same prompt. This technique is slightly slower than the first one, as it requires more function evaluations.

Image generation
Input Text→ Output Image
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